JoePenna/Dreambooth-Stable-Diffusion

Add support for Stable Diffusion 2.0 models

0xdevalias opened this issue · 14 comments

0xdevalias commented

rudimentary support for stable diffusion 2.0

MrCheeze/stable-diffusion-webui@069591b

Originally posted by @152334H in AUTOMATIC1111/stable-diffusion-webui#5011 (comment)

0xdevalias commented

https://github.com/hafriedlander/diffusers/blob/stable_diffusion_2/scripts/convert_original_stable_diffusion_to_diffusers.py

Notes:

  • Only tested on the two txt2img models, not inpaint / depth2img / upscaling
  • You will need to change your text embedding to use the penultimate layer too
  • It spits out a bunch of warnings about vision_model, but that's fine
  • I have no idea if this is right or not. It generates images, no guarantee beyond that. (Hence no PR - if you're patient, I'm sure the Diffusers team will do a better job than I have)

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

Here's an example of accessing the penultimate text embedding layer https://github.com/hafriedlander/stable-diffusion-grpcserver/blob/b34bb27cf30940f6a6a41f4b77c5b77bea11fd76/sdgrpcserver/pipeline/text_embedding/basic_text_embedding.py#L33

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

doesn't seem to work for me on the 768-v model using the v2 config for v

TypeError: EulerDiscreteScheduler.init() got an unexpected keyword argument 'prediction_type'

Originally posted by @devilismyfriend in huggingface/diffusers#1388 (comment)

You need to use the absolute latest Diffusers and merge this PR (or use my branch which has it in it) huggingface/diffusers#1386

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

(My branch is at https://github.com/hafriedlander/diffusers/tree/stable_diffusion_2)

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

0xdevalias commented

testing in progress on the horde https://github.com/Sygil-Dev/nataili/tree/v2
try it out Stable Diffusion 2.0 on our UI's

https://tinybots.net/artbot
https://aqualxx.github.io/stable-ui/
https://dbzer0.itch.io/lucid-creations

https://sigmoid.social/@stablehorde/109398715339480426

SD 2.0

  • Initial implementation ready for testing
  • img2img
  • inpainting
  • k_diffusers support

Originally posted by @AlRlC in https://github.com/Sygil-Dev/nataili/issues/67#issuecomment-1326385645

0xdevalias commented

Should work now, make sure you check the box "redownload original model" when choosing V2

https://colab.research.google.com/github/TheLastBen/fast-stable-diffusion/blob/main/fast_stable_diffusion_AUTOMATIC1111.ipynb

Requires more than 12GB of RAM for now, so free colab probably won't suffice.

Originally posted by @TheLastBen in TheLastBen/fast-stable-diffusion#599 (comment)

0xdevalias commented

From @pcuenca on the HF discord:

We are busy preparing a new release of diffusers to fully support Stable Diffusion 2. We are still ironing things out, but the basics already work from the main branch in github. Here's how to do it:

  • Install diffusers from github alongside its dependencies:
pip install --upgrade git+https://github.com/huggingface/diffusers.git transformers accelerate scipy
  • Use the code in this script to run your predictions:
from diffusers import DiffusionPipeline, EulerDiscreteScheduler
import torch

repo_id = "stabilityai/stable-diffusion-2"
device = "cuda"

scheduler = EulerDiscreteScheduler.from_pretrained(repo_id, subfolder="scheduler", prediction_type="v_prediction")
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16", scheduler=scheduler)
pipe = pipe.to(device)

prompt = "High quality photo of an astronaut riding a horse in space"
image = pipe(prompt, width=768, height=768, guidance_scale=9).images[0]
image.save("astronaut.png")

Originally posted by @vvvm23 in huggingface/diffusers#1392 (comment)

0xdevalias commented

how sure are you that your conversion is correct? I'm trying to diagnose a difference I get between your 768 weights and my conversion script. There's a big difference, and in general I much prefer the results from my conversion. It seems specific to the unet - if I replace my unet with yours I get the same results.

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

OK, differential diagnostic done, it's the Tokenizer. How did you create the Tokenizer at https://huggingface.co/stabilityai/stable-diffusion-2/tree/main/tokenizer? I just built a Tokenizer using AutoTokenizer.from_pretrained("laion/CLIP-ViT-H-14-laion2B-s32B-b79K") - it seems to give much better results.

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

I've put "my" version of the Tokenizer at https://huggingface.co/halffried/sd2-laion-clipH14-tokenizer/tree/main. You can just replace the tokenizer in any pipeline to test it if you're interested.

Originally posted by @hafriedlander in huggingface/diffusers#1388 (comment)

leondelee commented

when will Dreambooth support sd2

  • 👍5
0xdevalias commented

diffusers==0.9.0 with Stable Diffusion 2 is live!

https://github.com/huggingface/diffusers/releases/tag/v0.9.0

Originally posted by @anton-l in huggingface/diffusers#1388 (comment)

0xdevalias commented

I've almost finished a proper implementation of Stable Diffusion 2.0 in Automatic1111, so that it runs locally and automatically updates everything and works on 4GB lowvram. It supports both 1.5 and 2.0 models and you can switch between models from the menu like normal.

So far the 512x512 base model, 512x512 inpainting model, and the 768x768 v-prediction model work properly. The upscaler model and depth models load correctly but don't work to generate images yet.
It gives an error trying to load old Textual Inversion embeddings with the new models, but that can't be helped.
And the PLMS sampling method isn't working.
I'll push it soon.

Originally posted by @CarlKenner in AUTOMATIC1111/stable-diffusion-webui#5011 (comment)

  • ❤️1
0xdevalias commented

when will Dreambooth support sd2

While it's not dreambooth, this repo seems to have support for finetuning SDv2:

0xdevalias commented

And looking at the huggingface/diffusers repo, there are a few issues that seem to imply people may be getting dreambooth things working with that (or at least trying to), eg.:

goffi-contrib commented

For the depth model, the https://github.com/epitaque/dreambooth_depth2img repos does it by generating depth map for every input images, it would be great to see this integrated in this repos.